DiffTransfer: A Person Portrait Style Transfer Method Based on Stable Diffusion

Authors

  • Jizheng Wang

DOI:

https://doi.org/10.54097/4113h368

Keywords:

Style Transfer, diffusion model, multimodal.

Abstract

With the widespread adoption of high-performance GPUs, NPUs, and other devices in recent years, the threshold for deploying large models on the edge has gradually lowered. Meanwhile, as social networks continue to thrive and virtual reality technology advances rapidly, there is an urgent demand for individuals to possess a realistic virtual identity to meet the ever-evolving social needs. In this context, portrait style transfer has attracted growing research interest. Different from the previous works that mainly follow the pipeline of generative adversarial network or variational autoencoder, this research proposes a method for person portrait style transfer based on the stable diffusion model. By pre-training underlying models, users only need to upload a small number of photos to generate corresponding LORA models for individuals. The approach involves style fusion using textual prompts and other models to modify portrait styles and generate images of individuals in different poses. Extensive visualization results validate the effectiveness of this work, which is able to generate realistic portrait images of people with different styles.

Downloads

Download data is not yet available.

References

Portilla, J., Simoncelli, E.P. A Parametric Texture Model Based on Joint Statistics of Complex Wavelet Coefficients. International Journal of Computer Vision 40, 49 – 70 (2000).

Gatys, Leon A., Alexander S. Ecker, and Matthias Bethge. "A neural algorithm of artistic style." arXiv preprint arXiv: 1508.06576 (2015).

Ghiasi, Golnaz, et al. "Exploring the structure of a real-time, arbitrary neural artistic stylization network." arXiv preprint arXiv: 1705. 06830 (2017).

Chen, Shoufa, et al. "Diffusiondet: Diffusion model for object detection." Proceedings of the IEEE/CVF International Conference on Computer Vision. 2023.

Baranchuk, Dmitry, et al. "Label-efficient semantic segmentation with diffusion models." arXiv preprint arXiv: 2112. 03126 (2021).

Rombach, Robin, et al. "High-resolution image synthesis with latent diffusion models." Proceedings of the IEEE/CVF conference on computer vision and pattern recognition. 2022.

Andrew. How does Stable Diffusion work? Retrieved on November 18, 2023. Retrieved from https: //stable-diffusion-art.com/how-stable-diffusion-work/.

Radford, Alec, et al. "Learning transferable visual models from natural language supervision." International conference on machine learning. PMLR, 2021.

Hu, Edward J., et al. "Lora: Low-rank adaptation of large language models." arXiv preprint arXiv: 2106. 09685 (2021).

Ruiz, Nataniel, et al. "Dreambooth: Fine tuning text-to-image diffusion models for subject-driven generation." Proceedings of the IEEE/CVF Conference on Computer Vision and Pattern Recognition. 2023.

Gal, Rinon, et al. "An image is worth one word: Personalizing text-to-image generation using textual inversion." arXiv preprint arXiv: 2208. 01618 (2022).

Defazio, Aaron, and Konstantin Mishchenko. "Learning-rate-free learning by D-adaptation." arXiv preprint arXiv: 2301. 07733 (2023).

Zhang, Lvmin, Anyi Rao, and Maneesh Agrawala. "Adding conditional control to text-to-image diffusion models." Proceedings of the IEEE/CVF International Conference on Computer Vision. 2023.

Downloads

Published

13-03-2024